Our Discord : / discord . In this video, I am explaining how to extract the taught subject information from your trained model and inject into any new custom Stable Diffusion model. If I have been of assistance to you and you would like to show your support for my work, please consider becoming a patron on 🥰 / secourses
Playlist of Stable Diffusion Tutorials, Automatic1111 and Google Colab Guides, DreamBooth, Textual Inversion / Embedding, LoRA, AI Upscaling, Pix2Pix, Img2Img:
• Stable Diffusion Tutorials, Automatic...
Easiest Way to Install & Run Stable Diffusion Automatic1111 Web UI on PC by Using Open Source Automatic Installer
• Easiest Way to Install & Run Stable D...
How to use Stable Diffusion V2.1 and Different Models in the Automatic1111 Web UI - SD 1.5 vs 2.1 vs Anything V3
• How to use Stable Diffusion V2.1 and ...
Zero To Hero Stable Diffusion DreamBooth Tutorial By Using Automatic1111 Web UI - Ultra Detailed
• Zero To Hero Stable Diffusion DreamBo...
How To Do Stable Diffusion LORA Training By Using Automatic1111 Web UI On Different Models - Tested SD 1.5, SD 2.1
• How To Do Stable Diffusion LORA Train...
How to Do DreamBooth training on a Google Colab without needing any GPU and download trained ckpt model file
• Transform Your Selfie into a Stunning...
How to Use SD 2.1 & Custom Models on Google Colab for Training with Dreambooth & Image Generation
• How to Use SD 2.1 & Custom Models on ...
0:00 Introduction to how to inject / merge / combine your models by using checkpoint merger
1:48 Start of the tutorial
1:57 My face trained model used training dataset
2:12 The image quality of the default trained model (SD 1.5 official version)
2:44 How to inject your trained info from your trained model into a new custom model
3:04 What are primary model, secondary model and tertiary model
3:32 What is the strategy for extracting your trained subject from trained model and inject into a new custom model
4:31 What is Checkpoint Merger multiplier
5:01 Add Difference selection
5:25 How to use newly merged model
5:54 How to select proper prompt strength and CFG value for the new subject injected model
9:22 How to join our discord channel to ask anything and get support for free
Stable Diffusion is a powerful deep learning model that has the ability to generate highly detailed images from text descriptions. The model was developed by the CompVis group at LMU Munich, and was released in 2022 by a collaboration of several organizations, including Stability AI, CompVis LMU, and Runway, with support from EleutherAI and LAION.
One of the key features of Stable Diffusion is that it is a latent diffusion model, a type of deep generative neural network. This allows the model to generate a wide range of images from a given text prompt, and to make use of a large amount of data in order to improve its performance. The model is also highly efficient, and can run on most consumer hardware that is equipped with a moderate GPU with at least 8 GB VRAM, unlike many previous proprietary text-to-image models such as DALL-E and Midjourney which could only be accessed through cloud services.
In addition to its impressive capabilities, Stable Diffusion is also highly accessible to users. The model's code and model weights have been made publicly available, meaning that anyone can use the model and experiment with its capabilities. This opens up a wide range of possibilities for researchers and developers, who can use the model to explore new applications and improve its performance.
Another important feature of #StableDiffusion is the #DreamBooth approach which is a new way for personalizing text-to-image models. This method allows users to fine-tune a pre-trained text-to-image model like Imagen with just a few images of a specific subject. By doing this, the model learns to associate a unique identifier with that subject, allowing the generation of fully new and photorealistic images of the subject in a variety of contexts, poses and lighting conditions, while preserving the subject's unique features. This opens up a wide range of possibilities for the field, including text-guided view synthesis, appearance modification, and artistic rendering.
In October 2022, Stability AI raised US$101 million in a round led by Lightspeed Venture Partners and Coatue Management, which is a testament of how promising the technology is for investors. The Stable Diffusion is a powerful new tool that is poised to revolutionize the field of deep learning and text-to-image generation.
#Midjourney is a text-to-image generative model developed by OpenAI, which is known for its impressive capabilities in generating detailed images from text prompts. It is a variant of the DALL-E model which is a deep learning model that was trained on a massive dataset of images and captions, allowing it to learn the relationships text image data.